# Art Style Generation

This is a LoRA adapter trained for the FLUX.1-dev text-to-image model, compatible with diffusers or ComfyUI.

This is a LoRA adapter specifically designed for the FLUX.1-dev text-to-image generation model, compatible with diffusers or ComfyUI

A diffusion transformer model trained on the Wikiart dataset for generating artwork images

A diffusion transformer model trained on the Wikiart dataset for generating artistic style images

This is a fixed version of the Variational Autoencoder (VAE) for the Stable Diffusion XL (SDXL) model, focusing on FP16 precision optimization.

FLUX.1-dev is a text-to-image generation model based on Stable Diffusion technology, supporting LoRA fine-tuning, suitable for creative image generation tasks.

Lexica Aperture V3.5 is a free text-to-image generation model, similar to the paid version of Lexica Aperture V3.5, supporting multilingual input.

A text-to-image generation model based on stable diffusion technology, capable of producing high-quality images from textual descriptions.

A ControlNet checkpoint trained based on runwayml/stable-diffusion-v1-5, using fused color and edge detection as conditional input for image generation control.

This is a Dreambooth fine-tuned model based on Stable Diffusion v1.5, trained with 21 images for generating style-specific artwork.

ProllyMix is a text-to-image generation model based on Stable Diffusion, suitable for creative image generation tasks.

This is a diffusion model initialized with Google's DDPM bedroom image model and fine-tuned on the WikiArt dataset for unconditional image generation.

A stable diffusion-based text-to-image generation model supporting creative image generation